What problem does it solve?
This Skill helps you reliably turn text prompts (and optional reference images) into high-quality images using Stable Diffusion through HuggingFace Diffusers, removing the need to manually wire up complex diffusion pipelines.
Core Features & Use Cases
- Text-to-Image generation: Create images directly from natural-language prompts with controllable quality via steps and guidance scale.
- Image-to-Image translation: Transform an existing image using a prompt and a strength parameter to control how much to change.
- Inpainting: Fill masked regions with prompt-guided content for edits like removing objects or correcting parts of an image.
- Structured conditioning (ControlNet): Add spatial control (e.g., edges, pose, depth) to keep composition aligned with a reference.
- Performance and quality controls: Switch schedulers, enable memory optimizations (CPU offload, slicing, xFormers), and select model variants (SD 1.x/SDXL/SD 3.0/Flux).
- Quality and iteration workflows: Use SDXL refiner-style multi-stage generation patterns and reproducible seeds.
Quick Start
Use the stable-diffusion-image-generation skill to generate an image from the attached prompt by running text-to-image generation with Stable Diffusion via Diffusers.