What problem does it solve? Generating custom images programmatically requires navigating complex diffusion model pipelines, schedulers, memory constraints, and model variants. This Skill provides complete guidance for text-to-image generation, image-to-image translation, inpainting, and custom diffusion workflows using the HuggingFace Diffusers library. ## Core Features & Use Cases - Text-to-Image Generation: Create images from natural language prompts using SD 1.5, SDXL, SD 3.0, or Flux models with tunable guidance scale, steps, and negative prompts. - Image Transformation & Inpainting: Perform image-to-image translation, fill masked regions, and apply ControlNet spatial conditioning (edges, poses, depth) for precise control. - LoRA & Memory Optimization: Load LoRA style adapters, enable CPU offloading, attention slicing, and xFormers to run large models on limited GPU memory. - Use Case: A developer building a product mockup tool uses SDXL with a DPM scheduler and negative prompts to generate 1024x1024 marketing visuals from text descriptions, then deploys the pipeline behind a FastAPI endpoint. ## Quick Start Ask the AI to generate an image from a text prompt using Stable Diffusion with the Diffusers library, such as creating a 1024x1024 landscape with SDXL and a fast scheduler.